Local Prompt Adaptation for Style-Consistent Multi-Object Generation in Diffusion Models
On the map: Diffusion models → Image generation, text-to-image and editing
Abstract (arXiv)
Diffusion models have become a powerful backbone for text-to-image generation, producing high-quality visuals from natural language prompts. However, when prompts involve multiple objects alongside global or local style instructions, the outputs often drift in style and lose spatial coherence, limiting their reliability for controlled, style-consistent scene generation. We present Local Prompt Adaptation (LPA), a lightweight, training-free method that splits the prompt into content and style tokens, then injects them selectively into the U-Net's attention layers at chosen timesteps. By conditioning object tokens early and style tokens later in the denoising process, LPA improves both layout control and stylistic uniformity without additional training cost. We conduct extensive ablations across parser settings and injection windows, finding that the best configuration -- lpa late only with a 300-650 step window -- delivers the strongest balance of prompt alignment and style consistency. On the T2I benchmark, LPA improves CLIP-prompt alignment over vanilla SDXL by +0.41% and over SD1.5 by +0.34%, with no diversity loss. On our custom 50-prompt style-rich benchmark, LPA achieves +0.09% CLIP-prompt and +0.08% CLIP-style gains over baseline. Our method is model-agnostic, easy to integrate, and requires only a single configuration change, making it a practical choice for controllable, style-consistent multi-object generation.
Key facts. No citing papers in the map yet (data through 14 Oct 2026). Placed in Image generation, text-to-image and editing, one of 14 directions in Diffusion models. Builds on 37 of its 57 references that are in the map. First posted on arXiv 27 Jul 2025.
Builds on
37 of 57 references.
Closest ancestors · 10
Same direction, newest first
- 2023DragDiffusion: Harnessing Diffusion Models for Interactive Point-based Image Editing26 Jun
- 2023MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation16 Feb
- 2023Adding Conditional Control to Text-to-Image Diffusion Models10 Feb
- 2023Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models31 Jan
- 2022SpaText: Spatio-Textual Representation for Controllable Image Generation25 Nov
- 2022eDiff-I: Text-to-Image Diffusion Models with an Ensemble of Expert Denoisers2 Nov
- 2022Blended Latent Diffusion6 Jun
- 2022Compositional Visual Generation with Composable Diffusion Models3 Jun
- 2022Hierarchical Text-Conditional Image Generation with CLIP Latents13 Apr
- 2021GLIDE: Towards Photorealistic Image Generation and Editing with Text-Guided Diffusion Models20 Dec
Other references · 27
- 2017Attention Is All You Needcited by 60,361
- 2021Learning Transferable Visual Models From Natural Language Supervisioncited by 33,744
- 2020Denoising Diffusion Probabilistic Modelscited by 21,076
- 2021High-Resolution Image Synthesis with Latent Diffusion Modelscited by 18,863
- 2021LoRA: Low-Rank Adaptation of Large Language Modelscited by 14,380
- 2020Score-Based Generative Modeling through Stochastic Differential Equationscited by 9,122
- 2021Diffusion Models Beat GANs on Image Synthesiscited by 8,255
- 2015Deep Unsupervised Learning using Nonequilibrium Thermodynamicscited by 7,058
- 2022Scalable Diffusion Models with Transformerscited by 6,054
- 2023SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesiscited by 4,356
- 2021Improved Denoising Diffusion Probabilistic Modelscited by 3,761
- 2020Taming Transformers for High-Resolution Image Synthesiscited by 3,253
15 more
- 2018Large Scale GAN Training for High Fidelity Natural Image Synthesiscited by 2,626
- 2022Prompt-to-Prompt Image Editing with Cross Attention Controlcited by 2,431
- 2022An Image is Worth One Word: Personalizing Text-to-Image Generation using Textual Inversioncited by 2,245
- 2021SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equationscited by 1,937
- 2022RePaint: Inpainting using Denoising Diffusion Probabilistic Modelscited by 1,430
- 2021Grounded Language-Image Pre-trainingcited by 1,156
- 2021StyleCLIP: Text-Driven Manipulation of StyleGAN Imagerycited by 991
- 2019Improved Precision and Recall Metric for Assessing Generative Modelscited by 944
- 2018Image Transformercited by 772
- 2023Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generatorscited by 670
- 2021Vector-quantized Image Modeling with Improved VQGANcited by 658
- 2022DiffuSeq: Sequence to Sequence Text Generation with Diffusion Modelscited by 450
- 2022MCVD: Masked Conditional Video Diffusion for Prediction, Generation, and Interpolationcited by 309
- 2023Dreamix: Video Diffusion Models are General Video Editorscited by 191
- 2023Diffusion Models as Masked Autoencoderscited by 71
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- 2024Words Worth a Thousand Pictures: Measuring and Understanding Perceptual Variability in Text-to-Image Generation
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