Fine-Tuning Stable Diffusion XL for Stylistic Icon Generation: A Comparison of Caption Size
On the map: Diffusion models → Image generation, text-to-image and editing
Abstract (arXiv)
In this paper, we show different fine-tuning methods for Stable Diffusion XL; this includes inference steps, and caption customization for each image to align with generating images in the style of a commercial 2D icon training set. We also show how important it is to properly define what "high-quality" really is especially for a commercial-use environment. As generative AI models continue to gain widespread acceptance and usage, there emerge many different ways to optimize and evaluate them for various applications. Specifically text-to-image models, such as Stable Diffusion XL and DALL-E 3 require distinct evaluation practices to effectively generate high-quality icons according to a specific style. Although some images that are generated based on a certain style may have a lower FID score (better), we show how this is not absolute in and of itself even for rasterized icons. While FID scores reflect the similarity of generated images to the overall training set, CLIP scores measure the alignment between generated images and their textual descriptions. We show how FID scores miss significant aspects, such as the minority of pixel differences that matter most in an icon, while CLIP scores result in misjudging the quality of icons. The CLIP model's understanding of "similarity" is shaped by its own training data; which does not account for feature variation in our style of choice. Our findings highlight the need for specialized evaluation metrics and fine-tuning approaches when generating high-quality commercial icons, potentially leading to more effective and tailored applications of text-to-image models in professional design contexts.
Key facts. No citing papers in the map yet (data through 7 Oct 2026). Placed in Image generation, text-to-image and editing, one of 14 directions in Diffusion models. Builds on 11 of its 15 references that are in the map. First posted on arXiv 11 Jul 2024.
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- 2020Denoising Diffusion Probabilistic Modelscited by 20,851
- 2021High-Resolution Image Synthesis with Latent Diffusion Modelscited by 18,694
- 2023SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesiscited by 4,290
- 2013Auto-Encoding Variational Bayescited by 3,980
- 2021SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equationscited by 1,903
- 2023Rethinking FID: Towards a Better Evaluation Metric for Image Generationcited by 303
- 2015A note on the evaluation of generative modelscited by 291
- 2017GANs Trained by a Two Time-Scale Update Rule Converge to a Local Nash Equilibriumcited by 208
- 2023StyleDrop: Text-to-Image Generation in Any Stylecited by 197
- 2024BitsFusion: 1.99 bits Weight Quantization of Diffusion Modelcited by 37
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- 2023Navigating Text-To-Image Customization: From LyCORIS Fine-Tuning to Model Evaluation
- 2022The Stable Artist: Steering Semantics in Diffusion Latent Space
- 2024Direct Consistency Optimization for Robust Customization of Text-to-Image Diffusion Models
- 2024ArtWeaver: Advanced Dynamic Style Integration via Diffusion Model
- 2024Improving Text Generation on Images with Synthetic Captions
- 2023Improving Compositional Text-to-image Generation with Large Vision-Language Models
- 2025Training-free Stylized Text-to-Image Generation with Fast Inference
- 2024All Seeds Are Not Equal: Enhancing Compositional Text-to-Image Generation with Reliable Random Seeds