Morphological Addressing of Identity Basins in Text-to-Image Diffusion Models
On the map: Diffusion models → Image generation, text-to-image and editing
Abstract (arXiv)
We demonstrate that morphological pressure creates navigable gradients at multiple levels of the text-to-image generative pipeline. In Study~1, identity basins in Stable Diffusion 1.5 can be navigated using morphological descriptors -- constituent features like platinum blonde,'' beauty mark,'' and 1950s glamour'' -- without the target's name or photographs. A self-distillation loop (generating synthetic images from descriptor prompts, then training a LoRA on those outputs) achieves consistent convergence toward a specific identity as measured by ArcFace similarity. The trained LoRA creates a local coordinate system shaping not only the target identity but also its inverse: maximal away-conditioning produces eldritch'' structural breakdown in base SD1.5, while the LoRA-equipped model produces ``uncanny valley'' outputs -- coherent but precisely wrong. In Study~2, we extend this to prompt-level morphology. Drawing on phonestheme theory, we generate 200 novel nonsense words from English sound-symbolic clusters (e.g., cr-, sn-, -oid, -ax) and find that phonestheme-bearing candidates produce significantly more visually coherent outputs than random controls (mean Purity@1 = 0.371 vs. 0.209, p<0.00001p < 0.00001 p<0.00001, Cohen's d=0.55d = 0.55 d=0.55). Three candidates -- snudgeoid, crashax, and broomix -- achieve perfect visual consistency (Purity@1 = 1.0) with zero training data contamination, each generating a distinct, coherent visual identity from phonesthetic structure alone. Together, these studies establish that morphological structure -- whether in feature descriptors or prompt-level phonological form -- creates systematic navigational gradients through diffusion model latent spaces. We document phase transitions in identity basins, CFG-invariant identity stability, and novel visual concepts emerging from sub-lexical sound patterns.
Key facts. No citing papers in the map yet (data through 7 Oct 2026). Placed in Image generation, text-to-image and editing, one of 14 directions in Diffusion models. Builds on 9 of its 15 references that are in the map. First posted on arXiv 20 Feb 2026.
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- 2023Concept Sliders: LoRA Adaptors for Precise Control in Diffusion Models20 Nov
- 2023Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models31 Jan
- 2022Diffusion Art or Digital Forgery? Investigating Data Replication in Diffusion Models7 Dec
- 2022DreamBooth: Fine Tuning Text-to-Image Diffusion Models for Subject-Driven Generation25 Aug
Other references · 5
- 2021LoRA: Low-Rank Adaptation of Large Language Modelscited by 14,084
- 2022An Image is Worth One Word: Personalizing Text-to-Image Generation using Textual Inversioncited by 2,231
- 2023Extracting Training Data from Diffusion Modelscited by 723
- 2022Diffusion Models already have a Semantic Latent Spacecited by 324
- 2022DALLE-2 is Seeing Double: Flaws in Word-to-Concept Mapping in Text2Image Modelscited by 50
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